# Diffusion model optimization

LoRA fine-tuned model based on FLUX.1-dev, specifically designed for generating Ghibli-style images

A LoRA-based text-to-image diffusion model capable of generating Ghibli-style images

LoRA trained based on Wan2.1 14B T2V model, specifically designed for generating cinematic-grade fantasy landscape videos

Ghibli-style image generation LoRA model based on FLUX.1-dev base model

A LoRA fine-tuned text-to-image generation model based on FLUX.1-dev, specializing in generating specific images of the female protagonist Belle.

This is a LoRA model fine-tuned based on the FLUX.1-dev base model, used for generating images with specific styles.

LoRA model trained on the FLUX.1-dev base model for text-to-image generation, particularly suitable for stylized character scenes

A text-to-image generation model based on LoRA and diffusion models, specifically designed for generating realistic claymation-style images.

A text-to-image diffusion model based on LoRA technology, capable of generating hyper-realistic HD photos of Michelle Trachtenberg

A text-to-image diffusion model based on LoRA technology, capable of generating highly realistic and high-quality KayaS-style character images.

LoRA adaptation weights based on Bingsu/my-korean-stable-diffusion-v1-5, specifically fine-tuned for hanbok (Korean traditional clothing) image generation

A LoRA model fine-tuned based on the FLUX.1-schnell foundation model, specialized in generating cute comic-style images

A text-to-image generation model based on diffusion models and LoRA technology, capable of producing surreal-style photographs.

A text-to-image diffusion model based on LoRA technology, combining the features of DALL-E 3 and FLUX.1-dev for generating high-quality images.

DepthMaster is a refined single-step diffusion model that customizes generative features from diffusion models for discriminative depth estimation tasks.

A text-to-image generation model based on the FLUX.1-dev foundational model, specializing in line art style, particularly suitable for generating coloring book-like images.

A LoRA model based on FLUX.1-dev, specifically designed for generating 3D-style emoji images.

A LoRA model based on FLUX.1-dev, specializing in generating high-quality portrait images with fine control over various character features and backgrounds.

A LoRA adapter based on the FLUX.1-dev foundation model, specializing in generating high-quality 3D realistic Christmas-themed images.

CleanDIFT is a novel diffusion model feature extraction method that extracts noise-free and time-step-independent features by directly processing clean input images.

A LoRA adaptation model based on FLUX.1-dev foundation model, specialized for generating high-quality LEGO-themed images

Stable Diffusion 3.5 Large is a powerful text-to-image generation model based on diffusion model architecture, capable of producing high-quality images from text descriptions.

A LoRA model focused on generating amateur photography style images, suitable for diffusion model frameworks

This is a text-to-image diffusion model based on LoRA technology, specifically designed for generating cartoon-style images.

LoRA fine-tuned version based on the FLUX.1-dev foundation model, specializing in anime-style image generation

LoRA fine-tuned model based on FLUX.1-dev, specializing in anime-style image generation

A hyper-realistic model utilizing LoRA fine-tuning technology, capable of generating high-quality images with realistic textures, lighting, and intricate details.

A text-to-video prompt generation model fine-tuned based on Mistral-7B architecture, capable of generating complete prompts from minimal input words

An anime-style image generation model fine-tuned on Stable Diffusion, specializing in representing elements of Japanese anime, manga, and gaming culture.